Hello world! Welcome to my Stable Diffusion guide. The idea of this page is to guide people of every level into using my models and LoRA in the most effective way.
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In this section you'll find all the guides I've published divided by difficulty.
If you are a complete Stable Diffusion beginner, the first thing you'll need to do is to install the AUTOMATIC1111 UI. To do this i recommend you checking out my step-by-step tutorial on Medium:
Stable Diffusion Ultimate Guide pt. 1: Setup | by Umberto Grando | Mar, 2023 | Medium
In this tutorial you'll learn:
Once you've completed the setup, it's time to learn the secrets of prompting. The SD UI supports a series of commands that not everybody knows. Luckily for you I've written an article that explains everything step-by-step:
Stable Diffusion Ultimate Guide pt. 2: Prompting | by Umberto Grando | Mar, 2023 | Medium
In this article, you'll learn:
- Basic prompting: how to use a single prompt to generate text, and how to evaluate its quality and stability.
- Changing prompt weights: how to adjust the importance of each prompt keyword in relation to the others.
- Prompt Editing: how to change the number of steps that the model takes for a specific prompt.
- Prompt order: how to order your prompts according to their relevance and impact.
- Prompt matrix: how to use a matrix of prompts to test different variations in a single run.
Now that we know all the basics we need to learn how to create higher resolution images. SD is trained on 512x512 pictures, so trying to create big pictures usually results in poor quality or a lot of confusion. To do this we need to generate a smaller picture and then upscale it and re-generate:
Stable Diffusion Ultimate Guide pt. 3: High Resolution | by Umberto Grando | Mar, 2023 | Medium
In this tutorial you'll learn:
Sometimes you get an image that is 90% awesome and 10% horrible, that's when the inpainting tool come to the rescue. Check out my tutorial at:
Stable Diffusion Ultimate Guide pt. 4: Inpainting | by Umberto Grando | Apr, 2023 | Medium
In this tutorial you'll learn:
One of the problems you'll encounter when working with Stable Diffusion is the lack of controllability from the models. It's difficult to achieve a specific pose or add to a picture an object that SD doesn't know. Luckily with ControlNet you can achieve this easily.
You can check out my tutorial here:
Stable Diffusion Ultimate Guide pt. 5: Controlnet | by Umberto Grando | May, 2023 | Medium
In this tutorial you'll learn:
In this article you'll find my updated workflow to generate high-quality pictures using my custom models. Check out my tutorial at:
Stable Diffusion Ultimate Guide pt. 6: Workflow | by Umberto Grando | Jun, 2023 | Medium
In this tutorial you'll learn:
If your output looks like the image above you are not using a compatible VAE. You can install the one I'm currently using from here:
https://huggingface.co/iZELX1/Grapefruit/blob/main/Grapefruit.vae.pt
Once you've installed it (by saving it into the VAE folders of the UI), you need to go to the settings and under Stable Diffusion select the correct VAE.
If your pictures look less detailed when compared to mine the problem is probably the fact that you are not following a strong workflow. I've detailed mine in a step-by-step guide here:
Stable Diffusion Ultimate Guide pt. 3: High Resolution | by Umberto Grando | Mar, 2023 | Medium
Right now I don't have a specific workflow for creating LoRAs. I use the bmaltais/kohya_ss (github.com) GUI to create all my LoRAs.
The setting I (usually) use are:
If your image looks like the one above (with high saturation), you are probably using a high CFG value. To fix it you need to lower that value. The correct value changes from checkpoint to checkpoint, for example my Ruby model usually works with a value between 3 and 4 while my Anime model works well with a value of 6 or 7.
Models and LoRAs are listed on CivitAI:
Creator Profile | Civitai